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# Generated by project
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# IDEs
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.idea/
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.vscode/
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30
README.md
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README.md
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@ -1,12 +1,24 @@
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# Stable Diffusion 2.0
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# Stable Diffusion Version 2
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![t2i](assets/stable-samples/txt2img/768/merged-0006.png)
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![t2i](assets/stable-samples/txt2img/768/merged-0002.png)
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![t2i](assets/stable-samples/txt2img/768/merged-0005.png)
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This repository contains [Stable Diffusion](https://github.com/CompVis/stable-diffusion) models trained from scratch and will be continuously updated with
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new checkpoints. The following list provides an overview of all currently available models. More coming soon.
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## News
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**November 2022**
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**December 7, 2022**
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*Version 2.1*
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- New stable diffusion model (_Stable Diffusion 2.1-v_, [HuggingFace](https://huggingface.co/stabilityai/stable-diffusion-2-1)) at 768x768 resolution and (_Stable Diffusion 2.1-base_, [HuggingFace](https://huggingface.co/stabilityai/stable-diffusion-2-1-base)) at 512x512 resolution, both based on the same number of parameters and architecture as 2.0 and fine-tuned on 2.0, on a less restrictive NSFW filtering of the [LAION-5B](https://laion.ai/blog/laion-5b/) dataset.
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Per default, the attention operation of the model is evaluated at full precision when `xformers` is not installed. To enable fp16 (which can cause numerical instabilities with the vanilla attention module on the v2.1 model) , run your script with `ATTN_PRECISION=fp16 python <thescript.py>`
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**November 24, 2022**
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*Version 2.0*
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- New stable diffusion model (_Stable Diffusion 2.0-v_) at 768x768 resolution. Same number of parameters in the U-Net as 1.5, but uses [OpenCLIP-ViT/H](https://github.com/mlfoundations/open_clip) as the text encoder and is trained from scratch. _SD 2.0-v_ is a so-called [v-prediction](https://arxiv.org/abs/2202.00512) model.
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- The above model is finetuned from _SD 2.0-base_, which was trained as a standard noise-prediction model on 512x512 images and is also made available.
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- Added a [x4 upscaling latent text-guided diffusion model](#image-upscaling-with-stable-diffusion).
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@ -80,11 +92,11 @@ The weights are available via [the StabilityAI organization at Hugging Face](htt
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## Stable Diffusion v2.0
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## Stable Diffusion v2
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Stable Diffusion v2.0 refers to a specific configuration of the model
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Stable Diffusion v2 refers to a specific configuration of the model
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architecture that uses a downsampling-factor 8 autoencoder with an 865M UNet
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and OpenCLIP ViT-H/14 text encoder for the diffusion model. The _SD 2.0-v_ model produces 768x768 px outputs.
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and OpenCLIP ViT-H/14 text encoder for the diffusion model. The _SD 2-v_ model produces 768x768 px outputs.
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Evaluations with different classifier-free guidance scales (1.5, 2.0, 3.0, 4.0,
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5.0, 6.0, 7.0, 8.0) and 50 DDIM sampling steps show the relative improvements of the checkpoints:
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@ -97,16 +109,16 @@ Evaluations with different classifier-free guidance scales (1.5, 2.0, 3.0, 4.0,
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![txt2img-stable2](assets/stable-samples/txt2img/merged-0003.png)
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![txt2img-stable2](assets/stable-samples/txt2img/merged-0001.png)
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Stable Diffusion 2.0 is a latent diffusion model conditioned on the penultimate text embeddings of a CLIP ViT-H/14 text encoder.
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Stable Diffusion 2 is a latent diffusion model conditioned on the penultimate text embeddings of a CLIP ViT-H/14 text encoder.
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We provide a [reference script for sampling](#reference-sampling-script).
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#### Reference Sampling Script
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This script incorporates an [invisible watermarking](https://github.com/ShieldMnt/invisible-watermark) of the outputs, to help viewers [identify the images as machine-generated](scripts/tests/test_watermark.py).
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We provide the configs for the _SD2.0-v_ (768px) and _SD2.0-base_ (512px) model.
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We provide the configs for the _SD2-v_ (768px) and _SD2-base_ (512px) model.
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First, download the weights for [_SD2.0-v_](https://huggingface.co/stabilityai/stable-diffusion-2) and [_SD2.0-base_](https://huggingface.co/stabilityai/stable-diffusion-2-base).
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First, download the weights for [_SD2.1-v_](https://huggingface.co/stabilityai/stable-diffusion-2-1) and [_SD2.1-base_](https://huggingface.co/stabilityai/stable-diffusion-2-1-base).
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To sample from the _SD2.0-v_ model, run the following:
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To sample from the _SD2.1-v_ model, run the following:
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```
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python scripts/txt2img.py --prompt "a professional photograph of an astronaut riding a horse" --ckpt <path/to/768model.ckpt/> --config configs/stable-diffusion/v2-inference-v.yaml --H 768 --W 768
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@ -390,7 +390,7 @@ class DDPM(pl.LightningModule):
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elif self.parameterization == "v":
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target = self.get_v(x_start, noise, t)
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else:
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raise NotImplementedError(f"Paramterization {self.parameterization} not yet supported")
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raise NotImplementedError(f"Parameterization {self.parameterization} not yet supported")
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loss = self.get_loss(model_out, target, mean=False).mean(dim=[1, 2, 3])
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except:
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XFORMERS_IS_AVAILBLE = False
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# CrossAttn precision handling
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import os
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_ATTN_PRECISION = os.environ.get("ATTN_PRECISION", "fp32")
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def exists(val):
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return val is not None
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q, k, v = map(lambda t: rearrange(t, 'b n (h d) -> (b h) n d', h=h), (q, k, v))
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sim = einsum('b i d, b j d -> b i j', q, k) * self.scale
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# force cast to fp32 to avoid overflowing
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if _ATTN_PRECISION =="fp32":
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with torch.autocast(enabled=False, device_type = 'cuda'):
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q, k = q.float(), k.float()
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sim = einsum('b i d, b j d -> b i j', q, k) * self.scale
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else:
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sim = einsum('b i d, b j d -> b i j', q, k) * self.scale
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del q, k
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if exists(mask):
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10
modelcard.md
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modelcard.md
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@ -80,7 +80,7 @@ Stable Diffusion v2 mirrors and exacerbates biases to such a degree that viewer
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**Training Data**
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The model developers used the following dataset for training the model:
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- LAION-5B and subsets (details below). The training data is further filtered using LAION's NSFW detector, with a "p_unsafe" score of 0.1 (conservative). For more details, please refer to LAION-5B's [NeurIPS 2022](https://openreview.net/forum?id=M3Y74vmsMcY) paper and reviewer discussions on the topic.
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- LAION-5B and subsets (details below). The training data is further filtered using LAION's NSFW detector. For more details, please refer to LAION-5B's [NeurIPS 2022](https://openreview.net/forum?id=M3Y74vmsMcY) paper and reviewer discussions on the topic.
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**Training Procedure**
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Stable Diffusion v2 is a latent diffusion model which combines an autoencoder with a diffusion model that is trained in the latent space of the autoencoder. During training,
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- The output of the text encoder is fed into the UNet backbone of the latent diffusion model via cross-attention.
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- The loss is a reconstruction objective between the noise that was added to the latent and the prediction made by the UNet. We also use the so-called _v-objective_, see https://arxiv.org/abs/2202.00512.
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We currently provide the following checkpoints:
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We currently provide the following checkpoints, for various versions:
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### Version 2.1
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- `512-base-ema.ckpt`: Fine-tuned on `512-base-ema.ckpt` 2.0 with 220k extra steps taken, with `punsafe=0.98` on the same dataset.
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- `768-v-ema.ckpt`: Resumed from `768-v-ema.ckpt` 2.0 with an additional 55k steps on the same dataset (`punsafe=0.1`), and then fine-tuned for another 155k extra steps with `punsafe=0.98`.
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### Version 2.0
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- `512-base-ema.ckpt`: 550k steps at resolution `256x256` on a subset of [LAION-5B](https://laion.ai/blog/laion-5b/) filtered for explicit pornographic material, using the [LAION-NSFW classifier](https://github.com/LAION-AI/CLIP-based-NSFW-Detector) with `punsafe=0.1` and an [aesthetic score](https://github.com/christophschuhmann/improved-aesthetic-predictor) >= `4.5`.
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850k steps at resolution `512x512` on the same dataset with resolution `>= 512x512`.
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webdataset==0.2.5
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open-clip-torch==2.7.0
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gradio==3.11
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kornia==0.6
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invisible-watermark>=0.1.5
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streamlit-drawable-canvas==0.8.0
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-e .
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